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stable-diffusion-image-generation

7.6

by zechenzhangAGI

78Favorites
333Upvotes
0Downvotes

State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.

stable-diffusion

7.6

Rating

0

Installs

AI & LLM

Category

Quick Review

Exceptional skill documentation for Stable Diffusion image generation. The description clearly indicates when to use this skill (text-to-image, image-to-image, inpainting, custom pipelines), making it easily invokable by a CLI agent. Task knowledge is comprehensive with complete, runnable code examples covering basic generation, SDXL, schedulers, ControlNet, LoRA, memory optimization, and multiple workflows. Structure is excellent with clear sections, comparison tables, and appropriate delegation of advanced topics to reference files. The skill provides significant value by encapsulating complex diffusion model setup, scheduler selection, memory optimization strategies, and multi-component orchestration that would otherwise require extensive tokens and trial-and-error. Minor room for improvement in novelty score as some basic text-to-image tasks could be handled by simpler APIs, but the comprehensive coverage of advanced features (ControlNet, LoRA fusion, custom schedulers) and optimization techniques strongly justifies this as a dedicated skill.

LLM Signals

Description coverage9
Task knowledge10
Structure9
Novelty8

GitHub Signals

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Last commit 0 days ago

Publisher

zechenzhangAGI

zechenzhangAGI

Skill Author

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zechenzhangAGI

Skill Author

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