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stable-diffusion-image-generation

8.7

by davila7

127Favorites
446Upvotes
0Downvotes

State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.

stable-diffusion

8.7

Rating

0

Installs

AI & LLM

Category

Quick Review

Excellent comprehensive skill for Stable Diffusion image generation. The description clearly covers text-to-image, image-to-image, inpainting, and custom pipelines. Task knowledge is outstanding with complete working code examples for all major use cases, parameter explanations, optimization techniques, and multiple workflows. Structure is well-organized with clear sections, tables, and logical progression from basics to advanced topics, with additional details appropriately referenced in separate files. Novelty is strong as orchestrating Stable Diffusion pipelines with proper memory management, scheduler selection, ControlNet integration, and LoRA handling would require significant tokens and expertise from a CLI agent. Minor room for improvement: could explicitly mention hardware requirements upfront and provide a decision tree for model selection.

LLM Signals

Description coverage9
Task knowledge10
Structure9
Novelty8

GitHub Signals

18,073
1,635
132
71
Last commit 0 days ago

Publisher

davila7

davila7

Skill Author

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Publisher

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davila7

Skill Author

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